  Text-to-image diffusion models are a class of deep generative models that
have demonstrated an impressive capacity for high-quality image generation.
However, these models are susceptible to implicit biases that arise from
web-scale text-image training pairs and may inaccurately model aspects of
images we care about. This can result in suboptimal samples, model bias, and
images that do not align with human ethics and preferences. In this paper, we
present an effective scalable algorithm to improve diffusion models using
Reinforcement Learning (RL) across a diverse set of reward functions, such as
human preference, compositionality, and fairness over millions of images. We
illustrate how our approach substantially outperforms existing methods for
aligning diffusion models with human preferences. We further illustrate how
this substantially improves pretrained Stable Diffusion (SD) models, generating
samples that are preferred by humans 80.3% of the time over those from the base
SD model while simultaneously improving both the composition and diversity of
generated samples.
